In the rapidly evolving AI landscape, models for visual creation, such as Stable Diffusion, FLUX.1, HiDream-I1, ControlNet, Animagine XL, and Stable Video Diffusion, are transforming creative expression by allowing for the generation of photorealistic images, videos, and even anime-style visuals from text prompts. Stable Diffusion has become notably popular for its ability to generate images from text and image prompts using diffusion models, while FLUX.1, developed by former creators of Stable Diffusion, offers state-of-the-art performance in visual quality and prompt adherence. HiDream-I1 stands out for its ability to handle complex prompts and offers natural-language image editing capabilities. ControlNet enhances diffusion models by allowing precise control over image generation with minimal resource requirements. Animagine XL focuses on anime-style images, leveraging a tag-based prompting system for precision. Stable Video Diffusion provides open-source video generation, although it is still in the research phase. These models face challenges such as legal concerns over copyright, the need for computational resources, and the complexity of deploying them in production, but they also open up new possibilities for creative industries.