Deploying and using Stable Diffusion XL 1.0

Stable Diffusion XL 1.0 is a highly capable text-to-image model that can generate stunningly detailed and accurate images from simple prompts, requiring an A100 for invocation and taking around 8-12 seconds to produce an image. It can be deployed in two clicks from Baseten's model library or packaged as a Truss on GitHub. The model has the option of using the Stable Diffusion Refiner model to add accuracy to difficult-to-generate details like facial features and hands, which can be chosen with the use_refiner parameter.

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